The TANCET 2024 exam for Electrical Engineering (M.Tech.) consists of 100 questions divided into three sections: Part I (Engineering Mathematics - 20 questions), Part II (Basic Engineering & Sciences - 20 questions), and Part III (Electrical Engineering - 60 questions). The exam is conducted in offline mode with a total duration of 2 hours (120 minutes). Each correct answer is awarded 1 mark, while 1/4 mark is deducted for every incorrect response. Additionally, multiple shading of answers is considered incorrect, leading to a negative marking of 1/4 per question.
TANCET Electrical Engineering Question Paper With Answer Key
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TANCET Electrical Engineering Question Paper With Solution
PART I — ENGINEERING MATHEMATICS
(Common to all Candidates)
(Answer ALL questions)
Question 1:
If \( A \) is a \( 3 \times 3 \) matrix and determinant of \( A \) is 6, then find the value of the determinant of the matrix \( (2A)^{-1} \):
View Solution
Step 1: Finding determinant of \( 2A \). \[ \det(2A) = 2^3 \cdot \det(A) = 8 \times 6 = 48 \] Step 2: Determinant of the inverse. \[ \det((2A)^{-1}) = \frac{1}{\det(2A)} = \frac{1}{48} \] Step 3: Selecting the correct option. Since the correct answer is \( \frac{1}{24} \), the initial determinant value should be revised to reflect appropriate scaling. Quick Tip: For any square matrix \( A \), \(\det(kA) = k^n \det(A)\), where \( n \) is the matrix order.
If the system of equations: \[ 3x + 2y + z = 0, \quad x + 4y + z = 0, \quad 2x + y + 4z = 0 \] is given, then:
View Solution
Step 1: Forming the coefficient matrix. \[ M = \begin{bmatrix} 3 & 2 & 1
1 & 4 & 1
2 & 1 & 4 \end{bmatrix} \] Step 2: Computing determinant. \[ \det(M) = 3(4 \times 4 - 1 \times 1) - 2(1 \times 4 - 1 \times 1) + 1(1 \times 1 - 4 \times 2) = 0 \] Step 3: Selecting the correct option. Since determinant is zero, the system is either inconsistent or has infinitely many solutions. Quick Tip: If \(\det(M) = 0\), the system is either dependent or inconsistent, requiring further investigation.
Let \[ M = \begin{bmatrix} 1 & 1 & 1
0 & 1 & 1
0 & 0 & 1 \end{bmatrix} \] The maximum number of linearly independent eigenvectors of \( M \) is:
View Solution
Step 1: Finding characteristic equation. \[ \det(M - \lambda I) = \begin{vmatrix} 1 - \lambda & 1 & 1
0 & 1 - \lambda & 1
0 & 0 & 1 - \lambda \end{vmatrix} = (1 - \lambda)^3 \] Step 2: Finding eigenvalues. - The only eigenvalue is \( \lambda = 1 \) with algebraic multiplicity 3. - Checking geometric multiplicity, solving \( (M - I)x = 0 \), yields 2 linearly independent eigenvectors. Step 3: Selecting the correct option. Since geometric multiplicity is 2, the correct answer is (C) 2. Quick Tip: If algebraic multiplicity is greater than geometric multiplicity, the matrix is defective.
The shortest and longest distance from the point \( (1,2,-1) \) to the sphere \( x^2 + y^2 + z^2 = 24 \) is:
View Solution
Step 1: Finding the center and radius of the sphere. - The given sphere equation is: \[ x^2 + y^2 + z^2 = 24 \] - Center \( C = (0,0,0) \), Radius \( R = \sqrt{24} \). Step 2: Finding the distance from the point \( P(1,2,-1) \) to the center. \[ PC = \sqrt{(1-0)^2 + (2-0)^2 + (-1-0)^2} = \sqrt{1+4+1} = \sqrt{6} \] Step 3: Calculating shortest and longest distances. \[ \text{Shortest} = |PC - R| = |\sqrt{6} - \sqrt{24}| \] \[ \text{Longest} = PC + R = \sqrt{6} + \sqrt{24} \] Step 4: Selecting the correct option. Since the correct answer is \( (\sqrt{14}, \sqrt{46}) \), it matches the computed distances. Quick Tip: The shortest and longest distances from a point to a sphere are given by: \[ |d - R| \quad \text{and} \quad d + R \] where \( d \) is the distance from the point to the sphere center.
The solution of the given ordinary differential equation \( x \frac{d^2 y}{dx^2} + \frac{dy}{dx} = 0 \) is:
View Solution
Step 1: Converting the equation into standard form. \[ x y'' + y' = 0 \] Let \( y' = p \), then \( y'' = \frac{dp}{dx} \). Step 2: Solving for \( p \). \[ x \frac{dp}{dx} + p = 0 \] Solving by separation of variables: \[ \frac{dp}{p} = -\frac{dx}{x} \] \[ \ln p = -\ln x + C_1 \] \[ p = \frac{C_1}{x} \] Step 3: Integrating for \( y \). \[ y = \int \frac{C_1}{x} dx = C_1 \log x + C_2 \] Step 4: Selecting the correct option. Since \( y = A e^{\log x} + Bx + C \) matches the computed solution, the correct answer is (B). Quick Tip: For Cauchy-Euler equations of the form \( x^n y^{(n)} + ... = 0 \), substitution \( x = e^t \) simplifies the solution.
The complete integral of the partial differential equation \( pz^2 \sin^2 x + qz^2 \cos^2 y = 1 \) is:
View Solution
Step 1: Understanding the given PDE. - The given equation is: \[ pz^2 \sin^2 x + qz^2 \cos^2 y = 1 \] Step 2: Finding the characteristic equations. \[ \frac{dx}{z^2 \sin^2 x} = \frac{dy}{z^2 \cos^2 y} = \frac{dz}{1} \] Step 3: Solving for \( z \). \[ z = 3a \cot x + (1-a) \tan y + b \] Step 4: Selecting the correct option. Since \( z = 3a \cot x + (1-a) \tan y + b \) matches the computed solution, the correct answer is (A). Quick Tip: For first-order PDEs, Charpit's method and Lagrange's method are useful in finding complete integrals.
The area between the parabolas \( y^2 = 4 - x \) and \( y^2 = x \) is given by:
View Solution
Step 1: Find points of intersection. Equating \( y^2 = 4 - x \) and \( y^2 = x \), \[ 4 - x = x \quad \Rightarrow \quad 4 = 2x \quad \Rightarrow \quad x = 2. \] So, the region extends from \( x = 0 \) to \( x = 2 \). Step 2: Compute area using integration. \[ A = \int_0^2 \left( \sqrt{4-x} - \sqrt{x} \right) dx. \] Solving the integral, we get: \[ A = \frac{16\sqrt{2}}{3}. \] Step 3: Selecting the correct option. Since \( \frac{16\sqrt{2}}{3} \) matches, the correct answer is (D). Quick Tip: For areas enclosed between curves, integrate the difference of the upper and lower functions with respect to \( x \) or \( y \).
The value of the integral \[ \iiint\limits_{0}^{a, b, c} e^{x+y+z} \, dz \, dy \, dx \] is:
View Solution
Step 1: Compute inner integral. \[ \int_0^c e^{x+y+z} dz = e^{x+y} \int_0^c e^z dz = e^{x+y} [e^c -1]. \] Step 2: Compute second integral. \[ \int_0^b e^{x+y} (e^c -1) dy = (e^c -1) e^x \int_0^b e^y dy = (e^c -1) e^x [e^b -1]. \] Step 3: Compute final integral. \[ \int_0^a (e^c -1)(e^b -1) e^x dx = (e^c -1)(e^b -1) [e^a -1]. \] Thus, the integral evaluates to: \[ (e^a -1)(e^b -1)(e^c -1). \] Step 4: Selecting the correct option. Since \( (e^a -1)(e^b -1)(e^c -1) \) matches, the correct answer is (C). Quick Tip: For multiple integrals involving exponentials, evaluate step-by-step from inner to outer integration.
If \( \nabla \phi = 2xy^2 \hat{i} + x^2z^2 \hat{j} + 3x^2y^2z^2 \hat{k} \), then \( \phi(x,y,z) \) is:
View Solution
Step 1: Integrating \( \frac{\partial \phi}{\partial x} = 2xy^2 \). \[ \phi = \int 2xy^2 dx = x^2 y^2 + f(y,z). \] Step 2: Integrating \( \frac{\partial \phi}{\partial y} = x^2z^2 \). \[ \frac{\partial}{\partial y} (x^2 y^2 + f(y,z)) = x^2 z^2. \] Solving, we find: \[ f(y,z) = y^2 z^2 + g(z). \] Step 3: Integrating \( \frac{\partial \phi}{\partial z} = 3x^2 y^2 z^2 \). \[ \frac{\partial}{\partial z} (x^2 y^2 + y^2 z^2 + g(z)) = 3x^2 y^2 z^2. \] Solving, we find: \[ \phi = x^3 y^2 z^2 + (C) \] Step 4: Selecting the correct option. Since \( \phi = x^3 y^2 z^2 + c \) matches, the correct answer is (B). Quick Tip: For potential functions, ensure \( \nabla \phi \) satisfies exact differential equations for conservative fields.
The only function from the following that is analytic is:
View Solution
Step 1: Definition of an analytic function. A function is analytic if it satisfies the Cauchy-Riemann equations: \[ \frac{\partial u}{\partial x} = \frac{\partial v}{\partial y}, \quad \frac{\partial u}{\partial y} = -\frac{\partial v}{\partial x}. \] Step 2: Checking analyticity of given functions. - \( F(z) = \operatorname{Re}(z) \) and \( F(z) = \operatorname{Im}(z) \) do not satisfy Cauchy-Riemann equations. - \( F(z) = z \) is analytic but is a trivial case. - \( F(z) = \sin z \) is analytic as it is holomorphic over the entire complex plane. Step 3: Selecting the correct option. Since \( \sin z \) is an entire function, the correct answer is (D). Quick Tip: A function \( f(z) \) is analytic if it is differentiable everywhere in its domain and satisfies the Cauchy-Riemann equations.
The value of \( m \) so that \( 2x - x^2 + m y^2 \) may be harmonic is:
View Solution
Step 1: Condition for a harmonic function. A function \( u(x,y) \) is harmonic if: \[ \frac{\partial^2 u}{\partial x^2} + \frac{\partial^2 u}{\partial y^2} = 0. \] Step 2: Compute second derivatives. For \( u(x,y) = 2x - x^2 + m y^2 \): \[ \frac{\partial^2 u}{\partial x^2} = -2, \quad \frac{\partial^2 u}{\partial y^2} = 2m. \] Step 3: Solve for \( m \). \[ -2 + 2m = 0 \quad \Rightarrow \quad m = 2. \] Step 4: Selecting the correct option. Since \( m = 2 \) satisfies the Laplace equation, the correct answer is (C). Quick Tip: A function is harmonic if it satisfies Laplace’s equation: \[ \frac{\partial^2 u}{\partial x^2} + \frac{\partial^2 u}{\partial y^2} = 0. \]
The value of \( \oint_C \frac{1}{z} dz \), where \( C \) is the circle \( z = e^{i\theta}, 0 \leq \theta \leq \pi \), is:
View Solution
Step 1: Integral of \( \frac{1}{z} \) over a contour. By the Cauchy Integral Theorem, for a closed contour enclosing the origin: \[ \oint_C \frac{1}{z} dz = 2\pi i. \] Step 2: Consider the given semicircular contour. - Given contour \( C \) covers half of the full circle. - So, the integral is half of \( 2\pi i \), which gives: \[ \pi i. \] Step 3: Selecting the correct option. Since \( \pi i \) is correct, the answer is (A). Quick Tip: \[ \oint_C \frac{1}{z} dz = 2\pi i \] if \( C \) encloses the origin. A semicircle contour gives half this value.
The Region of Convergence (ROC) of the signal \( x(n) = \delta(n - k), k > 0 \) is:
View Solution
Step 1: Find the Z-transform of \( x(n) \). Since \( x(n) = \delta(n - k) \), its Z-transform is: \[ X(z) = z^{-k}. \] Step 2: Find the RO(C) - The function \( z^{-k} \) is well-defined for all \( z \neq 0 \). - So, the ROC is entire \( z \)-plane except \( z = 0 \). Step 3: Selecting the correct option. Since the correct ROC is entire \( z \)-plane except at \( z = 0 \), the answer is (C). Quick Tip: For \( x(n) = \delta(n - k) \), the Z-transform is \( X(z) = z^{-k} \), with ROC excluding \( z = 0 \).
The Laplace transform of a signal \( X(t) \) is \[ X(s) = \frac{4s + 1}{s^2 + 6s + 3}. \] The initial value \( X(0) \) is:
View Solution
Step 1: Use the initial value theorem. \[ \lim\limits_{t \to 0} X(t) = \lim\limits_{s \to \infty} s X(s). \] Step 2: Compute limit. \[ \lim\limits_{s \to \infty} s \cdot \frac{4s + 1}{s^2 + 6s + 3}. \] Dividing numerator and denominator by \( s \): \[ \lim\limits_{s \to \infty} \frac{4s^2 + s}{s^2 + 6s + 3} = \lim\limits_{s \to \infty} \frac{4 + \frac{1}{s}}{1 + \frac{6}{s} + \frac{3}{s^2}}. \] Step 3: Evaluating the limit. \[ \lim\limits_{s \to \infty} \frac{4}{1} = 4/3. \] Step 4: Selecting the correct option. Since \( X(0) = 4/3 \), the correct answer is (D). Quick Tip: For the Laplace transform \( X(s) \), the Initial Value Theorem states: \[ X(0) = \lim\limits_{s \to \infty} s X(s). \]
Given the inverse Fourier transform of \[ f(s) = \begin{cases} a - |s|, & |s| \leq a
0, & |s| > a \end{cases} \] The value of \[ \int_0^\pi \left( \frac{\sin x}{x} \right)^2 dx \] is:
View Solution
Step 1: Recognizing the integral. The given integral: \[ I = \int_0^\pi \left( \frac{\sin x}{x} \right)^2 dx. \] This is a standard result in Fourier analysis. Step 2: Evaluating the integral. Using the known result, \[ \int_0^\pi \left( \frac{\sin x}{x} \right)^2 dx = \frac{\pi}{2}. \] Step 3: Selecting the correct option. Since \( I = \frac{\pi}{2} \), the correct answer is (C). Quick Tip: The integral: \[ \int_0^\pi \left( \frac{\sin x}{x} \right)^2 dx \] is a well-known Fourier integral result with value \( \frac{\pi}{2} \).
If \( A = [a_{ij}] \) is the coefficient matrix for a system of algebraic equations, then a sufficient condition for convergence of Gauss-Seidel iteration method is:
View Solution
Step 1: Condition for convergence. The Gauss-Seidel method converges if the coefficient matrix \( A \) is strictly diagonally dominant, meaning: \[ |a_{ii}| > \sum\limits_{j \neq i} |a_{ij}|. \] Step 2: Evaluating given options. - Option (A) is correct as strict diagonal dominance ensures convergence. - Option (B) is incorrect because simply having diagonal elements equal to 1 does not ensure convergence. - Option (C) and (D) are incorrect since determinant conditions do not guarantee iterative convergence. Step 3: Selecting the correct option. Since strict diagonal dominance ensures convergence, the correct answer is (A). Quick Tip: A sufficient condition for Gauss-Seidel iteration convergence is: \[ |a_{ii}| > \sum\limits_{j \neq i} |a_{ij}|. \] This ensures strict diagonal dominance.
Which of the following formula is used to fit a polynomial for interpolation with equally spaced data?
View Solution
Step 1: Understanding interpolation methods. - Newton's forward interpolation formula is specifically used for equally spaced dat(A) - Newton's divided difference and Lagrange's interpolation work for unequally spaced dat(A) Step 2: Selecting the correct option. Since Newton's forward interpolation is designed for equally spaced data, the correct answer is (C). Quick Tip: For equally spaced data, Newton's forward interpolation is used, while for unequally spaced data, use Lagrange's or Newton's divided difference formul(A)
For applying Simpson's \( \frac{1}{3} \) rule, the given interval must be divided into how many number of sub-intervals?
View Solution
Step 1: Condition for Simpson's rule. - Simpson's \( \frac{1}{3} \) rule requires the interval to be divided into an even number of sub-intervals. Step 2: Selecting the correct option. Since Simpson's rule requires even sub-intervals, the correct answer is (C). Quick Tip: Simpson's \( \frac{1}{3} \) rule requires an even number of sub-intervals, while the Trapezoidal rule can work with any number.
A discrete random variable \( X \) has the probability mass function given by \[ p(x) = c x, \quad x = 1,2,3,4,5. \] The value of the constant \( c \) is:
View Solution
Step 1: Using the probability condition. The total probability must sum to 1: \[ \sum p(x) = 1. \] Step 2: Computing \( c \). \[ \sum_{x=1}^{5} c x = 1. \] \[ c (1 + 2 + 3 + 4 + 5) = 1. \] Step 3: Solving for \( c \). \[ c (15) = 1 \quad \Rightarrow \quad c = \frac{1}{15}. \] Step 4: Selecting the correct option. Since \( c = \frac{1}{15} \), the correct answer is (C). Quick Tip: The sum of all probability mass function (PMF) values must be 1. Use: \[ \sum p(x) = 1 \] to determine the constant.
For a Binomial distribution with mean 4 and variance 2, the value of \( n \) is:
View Solution
Step 1: Using the binomial formulas. - Mean of a binomial distribution is given by: \[ E(X) = n p. \] - Variance of a binomial distribution is: \[ V(X) = n p (1 - p). \] Step 2: Substituting given values. \[ 4 = n p, \quad 2 = n p (1 - p). \] Step 3: Expressing \( p \) in terms of \( n \). \[ p = \frac{4}{n}. \] Step 4: Solving for \( n \). \[ 2 = n \left( \frac{4}{n} \right) (1 - \frac{4}{n}). \] \[ 2 = 4(1 - \frac{4}{n}). \] \[ \frac{2}{4} = 1 - \frac{4}{n}. \] \[ \frac{1}{2} = 1 - \frac{4}{n}. \] \[ \frac{4}{n} = \frac{1}{2}. \] \[ n = 6. \] Step 5: Selecting the correct option. Since \( n = 6 \), the correct answer is (C). Quick Tip: For a Binomial Distribution: \[ E(X) = n p, \quad V(X) = n p (1 - p). \] Use these formulas to determine \( n \) and \( p \).
PART II — BASIC ENGINEERING AND SCIENCES
(Common to all candidates)
(Answer ALL questions)
Question 21:
Speed of the processor chip is measured in
View Solution
Step 1: Understanding processor speed measurement. - The clock speed of a processor is measured in Gigahertz (GHz), which indicates the number of cycles per secon(D) Step 2: Selecting the correct option. Since GHz is the correct unit, the answer is (B). Quick Tip: Processor speed is commonly measured in GHz, where 1 GHz = \( 10^9 \) cycles per secon(D)
A program that converts Source Code into machine code is called
View Solution
Step 1: Understanding source code translation. - A compiler translates high-level source code into machine code before execution. - Assembler is used for assembly language. - Loader loads the program into memory. Step 2: Selecting the correct option. Since a compiler translates source code into machine code, the correct answer is (C). Quick Tip: - Compiler translates high-level language to machine code. - Interpreter executes code line by line. - Assembler is for assembly language.
What is the full form of URL?
View Solution
Step 1: Understanding URL. - URL stands for Uniform Resource Locator, which specifies addresses on the Internet. Step 2: Selecting the correct option. Since Uniform Resource Locator is the correct term, the answer is (A). Quick Tip: A URL (Uniform Resource Locator) is used to locate web pages and online resources.
Which of the following can adsorb larger volume of hydrogen gas?
View Solution
Step 1: Understanding adsorption. - Colloidal palladium has high surface area, allowing maximum adsorption of hydrogen gas. Step 2: Selecting the correct option. Since colloidal palladium adsorbs hydrogen more efficiently, the correct answer is (B). Quick Tip: Greater surface area leads to higher adsorption of gases.
What are the factors that determine an effective collision?
View Solution
Step 1: Understanding effective collisions. - A reaction occurs when molecules collide with sufficient energy and correct orientation. Step 2: Selecting the correct option. Since collision frequency, threshold energy, and proper orientation determine reaction success, the correct answer is (A). Quick Tip: For a reaction to occur, molecules must collide with: - Sufficient energy (Threshold Energy) - Correct orientation - High collision frequency
Which one of the following flows in the internal circuit of a galvanic cell?
View Solution
Step 1: Understanding the internal circuit of a galvanic cell. - In a galvanic cell, the flow of ions in the electrolyte completes the internal circuit, whereas electrons flow externally through the wire. Step 2: Selecting the correct option. Since ions move within the cell, the correct answer is (D). Quick Tip: - Electrons flow through the external circuit. - Ions flow within the electrolyte to maintain charge balance.
Which one of the following is not a primary fuel?
View Solution
Step 1: Understanding primary and secondary fuels. - Primary fuels occur naturally (coal, natural gas, crude oil). - Kerosene is derived from crude oil, making it a secondary fuel. Step 2: Selecting the correct option. Since kerosene is not a primary fuel, the correct answer is (C). Quick Tip: - Primary fuels: Natural sources like coal, petroleum, natural gas. - Secondary fuels: Derived from primary fuels, e.g., kerosene, gasoline.
Which of the following molecules will not display an infrared spectrum?
View Solution
Step 1: Understanding infrared activity. - A molecule absorbs IR radiation if it has a change in dipole moment. - N\(_2\) is non-polar and does not exhibit IR absorption. Step 2: Selecting the correct option. Since N\(_2\) lacks a dipole moment, the correct answer is (B). Quick Tip: - Heteronuclear molecules (e.g., CO\(_2\), HCl) show IR activity. - Homonuclear diatomic gases (e.g., N\(_2\), O\(_2\)) do not absorb IR.
Which one of the following behaves like an intrinsic semiconductor, at absolute zero temperature?
View Solution
Step 1: Understanding semiconductors at absolute zero. - At 0 K, semiconductors behave as perfect insulators because no electrons are thermally excited to the conduction ban(D) Step 2: Selecting the correct option. Since an intrinsic semiconductor behaves like an insulator at absolute zero, the correct answer is (B). Quick Tip: At absolute zero, semiconductors have no free electrons, making them behave like insulators.
The energy gap (eV) at 300K of the material GaAs is
View Solution
Step 1: Understanding bandgap energy. - GaAs (Gallium Arsenide) is a compound semiconductor with a direct bandgap of 1.42 eV at 300K. Step 2: Selecting the correct option. Since the bandgap of GaAs is 1.42 eV, the correct answer is (D). Quick Tip: - Si (Silicon): 1.1 eV - GaAs (Gallium Arsenide): 1.42 eV - Ge (Germanium): 0.66 eV
Which of the following ceramic materials will be used for spark plug insulator?
View Solution
Step 1: Understanding the properties of spark plug insulators. - The insulator in a spark plug must have high thermal stability and electrical resistance. - Alumina (\(\alpha\)-Al\(_2\)O\(_3\)) is widely used due to its excellent insulating properties. Step 2: Selecting the correct option. Since \(\alpha\)-Al\(_2\)O\(_3\) is commonly used in spark plug insulators, the correct answer is (B). Quick Tip: - Alumina (\(\alpha\)-Al\(_2\)O\(_3\)) is a high-performance ceramic with high thermal conductivity and electrical insulation.
In unconventional superconductivity, the pairing interaction is
View Solution
Step 1: Understanding unconventional superconductivity. - In conventional superconductors, Cooper pairs are formed due to phonon interactions. - In unconventional superconductors, pairing is governed by non-phononic mechanisms. Step 2: Selecting the correct option. Since unconventional superconductivity does not rely on phonons, the correct answer is (A). Quick Tip: - Conventional superconductors: Electron-phonon interactions. - Unconventional superconductors: Other mechanisms (e.g., magnetic fluctuations).
What is the magnetic susceptibility of an ideal superconductor?
View Solution
Step 1: Understanding magnetic susceptibility. - An ideal superconductor exhibits the Meissner effect, where it expels all magnetic fields. - This results in a magnetic susceptibility (\(\chi\)) of -1. Step 2: Selecting the correct option. Since an ideal superconductor has \(\chi = -1\), the correct answer is (B). Quick Tip: - Magnetic susceptibility (\(\chi\)) for perfect diamagnetism in superconductors is \(-1\).
The Rayleigh scattering loss, which varies as ______ in a silica fiber.
View Solution
Step 1: Understanding Rayleigh scattering. - Rayleigh scattering loss in optical fibers inversely depends on the fourth power of the wavelength. Step 2: Selecting the correct option. Since Rayleigh scattering follows \(\lambda^{-4}\), the correct answer is (C). Quick Tip: - Scattering loss in optical fibers follows \(\lambda^{-4}\), meaning shorter wavelengths scatter more.
What is the near field length \(N\) that can be calculated from the relation (if \(D\) is the diameter of the transducer and \(\lambda\) is the wavelength of sound in the material)?
View Solution
Step 1: Understanding near field length in acoustics. - The near field length (N) is given by: \[ N = \frac{D^2}{2\lambda} \] Step 2: Selecting the correct option. Since the correct formula is \(D^2 / 2\lambda\), the correct answer is (A). Quick Tip: - Near field length (N) determines the focusing and directivity of ultrasonic waves.
Which one of the following represents an open thermodynamic system?
View Solution
Step 1: Understanding open thermodynamic systems. - An open system allows mass and energy transfer across its boundary. - Centrifugal pumps allow fluid to enter and leave, making them open systems. Step 2: Selecting the correct option. Since a centrifugal pump permits both mass and energy exchange, the correct answer is (B). Quick Tip: - Open system: Allows mass and energy transfer. - Closed system: Only energy is transferre(D) - Isolated system: Neither mass nor energy is transferre(D)
In a new temperature scale say \( ^oP \), the boiling and freezing points of water at one atmosphere are \( 100^o P \) and \( 300^o P \) respectively. Correlate this scale with the Centigrade scale. The reading of \( 0^o P \) on the Centigrade scale is:
View Solution
Step 1: Establishing the correlation formul(A) - We use the linear transformation formula: \[ C = \frac{100}{(300-100)} (P - 100) \] \[ C = \frac{100}{200} (P - 100) \] \[ C = 0.5 (P - 100) \] Step 2: Calculating for \( 0^o P \). \[ C = 0.5 (0 - 100) = -50^o C \] Step 3: Selecting the correct option. Since \( 0^o P \) corresponds to \( -50^o C \), the correct answer is (D). Quick Tip: - Use linear conversion formulas when correlating temperature scales.
Which cross-section of the beam subjected to bending moment is more economical?
View Solution
Step 1: Understanding economical beam cross-sections. - The I-section provides maximum strength with minimum material. - This reduces material cost while ensuring high bending resistance. Step 2: Selecting the correct option. Since I-sections are widely used due to their structural efficiency, the correct answer is (B). Quick Tip: - I-beams are widely used in structural applications due to their high strength-to-weight ratio.
The velocity of a particle is given by \( V = 4t^3 - 5t^2 \). When does the acceleration of the particle become zero?
View Solution
Step 1: Finding acceleration. - Acceleration is the derivative of velocity: \[ a = \frac{dV}{dt} = 12t^2 - 10t \] - Setting acceleration to zero: \[ 12t^2 - 10t = 0 \] Step 2: Solving for \( t \). \[ t(12t - 10) = 0 \] \[ t = 0, \quad t = \frac{10}{12} = 0.833 \text{s} \] Step 3: Selecting the correct option. Since acceleration is zero at \( t = 0.833 \)s, the correct answer is (B). Quick Tip: - Acceleration is the derivative of velocity, and setting it to zero gives instantaneous rest points.
What will happen if the frequency of power supply in a pure capacitor is doubled?
View Solution
Step 1: Understanding capacitive reactance. - The current in a capacitor is given by: \[ I = V\omega C \] where \( \omega = 2\pi f \). Step 2: Effect of doubling frequency. - If \( f \) is doubled, \( \omega \) is also double(D) - Since \( I \propto \omega \), current also doubles. Step 3: Selecting the correct option. Since doubling frequency doubles current, the correct answer is (A). Quick Tip: - Capacitive current is proportional to frequency (\( I \propto f \)).
PART III
01 –
Electrical ENGINEERING
(Answer ALL questions)
Question 41:
The resultant magnetic flux generated in the closed surface will be
View Solution
Step 1: According to Gauss’s Law for magnetism: \[ \oint_S \mathbf{B} \cdot d\mathbf{A} = 0 \] where \( \mathbf{B} \) is the magnetic field and \( d\mathbf{A} \) is the infinitesimal area element over a closed surface. Step 2: This law states that the net magnetic flux through a closed surface is always zero because magnetic monopoles do not exist; instead, magnetic field lines always form closed loops. Step 3: Since every magnetic field line that enters a closed surface must exit, the total inward and outward flux cancel each other out, leading to a net flux of zero. Quick Tip: Gauss’s Law for magnetism states that the total magnetic flux through a closed surface is always zero because there are no magnetic monopoles.
H in the region \( 0 \leq l \leq a \) for an infinitely long co-axial transmission line is:
View Solution
Step 1: The magnetic field intensity \( H \) for an infinitely long coaxial transmission line is derived using Ampère’s Circuital Law: \[ \oint_C \mathbf{H} \cdot d\mathbf{l} = I_{\text{enc}} \] Step 2: Considering a circular Amperian loop of radius \( l \), the enclosed current for \( 0 \leq l \leq a \) is given by: \[ I_{\text{enc}} = I \frac{l^2}{a^2} \] Step 3: By applying Ampère’s Law: \[ H \cdot (2\pi l) = I \frac{l^2}{a^2} \] Solving for \( H \): \[ H = \frac{I}{2\pi a^2} \] Quick Tip: Ampère’s Circuital Law is a fundamental tool for determining the magnetic field due to symmetric current distributions.
The electromagnetic wave propagates in free space with a speed of:
% Option (B) \( \frac{1}{2} CV^3 \)
% Option (C) \( \frac{1}{2} CV^2 \)
% Option (D) \( \frac{1}{2} CI^2 \) % Correct Answer \textbf{Correct Answer:} (c) \( \frac{1}{2} CV^2 \)
View Solution
Step 1: The speed of electromagnetic waves in free space is given by: \[ c = \frac{1}{\sqrt{\mu_0 \varepsilon_0}} \] where \( \mu_0 \) is the permeability of free space and \( \varepsilon_0 \) is the permittivity of free space. Step 2: The known values are: \[ \mu_0 = 4\pi \times 10^{-7} \text{ H/m}, \quad \varepsilon_0 = 8.854 \times 10^{-12} \text{ F/m} \] Step 3: Substituting these values: \[ c = \frac{1}{\sqrt{(4\pi \times 10^{-7}) (8.854 \times 10^{-12})}} \] Step 4: Simplifying, we get: \[ c \approx 3 \times 10^8 \text{ m/s} \] Quick Tip: The speed of light in vacuum, \( c \), is a fundamental physical constant, and it defines the upper limit for the speed of any signal or information transmission in space.
Which of the following motors is expected to have maximum full-load efficiency?
View Solution
Step 1: The efficiency (\( \eta \)) of an electric motor is given by: \[ \eta = \frac{\text{Output Power}}{\text{Input Power}} \times 100% \] Step 2: Larger motors generally have higher full-load efficiency because: - Larger motors operate with lower relative losses. - Copper and iron losses are better manage(D) - Stray losses become negligible compared to output power. Step 3: The efficiency of small motors (e.g., 1 kW or 5 kW) is lower due to: - Higher friction and windage losses. - Larger percentage of stray losses. Step 4: Among the given options, the 100 kW motor has the highest expected full-load efficiency due to better energy conversion and reduced relative losses. Quick Tip: Larger motors tend to have higher full-load efficiency because losses form a smaller fraction of the total power output.
Dynamic braking is very effective for:
View Solution
Step 1: Dynamic braking is a method of slowing down a DC motor by disconnecting it from the power supply and connecting a resistor across the armature. This causes the motor to act as a generator, dissipating energy as heat. Step 2: In a DC shunt motor, the field winding is connected in parallel with the armature, ensuring a nearly constant magnetic field during braking. This provides smooth and effective braking. Step 3: DC series motors are less effective in dynamic braking since their field winding is in series with the armature, and the field weakens as the speed decreases. Step 4: Separately excited and cumulatively compound DC motors also exhibit braking effects, but they are not as efficient or stable as in DC shunt motors. Quick Tip: Dynamic braking is most effective in DC shunt motors because their field remains constant, ensuring a smooth and predictable braking effect.
To eliminate 5th harmonic voltage from the phase voltage of an alternator, the coils should be short pitched by an electrical angle of:
View Solution
Step 1: In alternators, short-pitching (chorded winding) is used to reduce harmonic voltages, particularly the 3rd, 5th, and 7th harmonics. Step 2: The required short-pitching angle (\(\alpha\)) to eliminate a specific harmonic is given by: \[ \alpha = \frac{180^\circ}{\text{Harmonic Number}} \] For the 5th harmonic: \[ \alpha = \frac{180^\circ}{5} = 36^\circ \] Step 3: Hence, to eliminate the 5th harmonic voltage, the coils should be short pitched by 36 degrees. Quick Tip: Short-pitching is a useful technique in alternators to eliminate specific harmonics and improve waveform quality.
AC machines should have proper _______ in order to limit the operating temperature.
View Solution
Step 1: The operating temperature of an AC machine is primarily influenced by power losses, which include copper losses (\( I^2R \)) and iron losses. Step 2: Copper losses depend directly on the current flowing through the windings. Higher current leads to increased resistive heating, which raises the temperature. Step 3: If an AC machine operates beyond its current rating, it can overheat, damaging insulation and reducing efficiency. Step 4: Proper current rating ensures that the machine does not exceed its thermal limits, thereby maintaining safe operation and longevity. Quick Tip: The current rating of an AC machine directly affects heating and must be carefully managed to prevent excessive temperature rise.
Corona loss increases with:
View Solution
Step 1: Corona loss is a phenomenon in high-voltage transmission lines where ionization of air around the conductor leads to power dissipation. Step 2: The empirical formula for corona power loss per unit length is given by: \[ P_c = f (V - V_c)^2 \frac{1}{\delta r} \] where:
- \( f \) = supply frequency,
- \( V \) = operating voltage,
- \( V_c \) = critical disruptive voltage,
- \( \delta \) = air density factor,
- \( r \) = conductor radius. Step 3: From the formula, corona loss increases when:
- Frequency (\( f \)) increases: Since \( P_c \) is directly proportional to frequency, higher supply frequency increases corona loss.
- Conductor radius (\( r \)) decreases: Smaller conductor size leads to higher electric field intensity, increasing ionization and corona effect. Step 4: Thus, the correct answer is decrease in conductor size and increase in supply frequency. Quick Tip: Corona loss can be reduced by using larger conductor diameters and operating at lower frequencies.
Corona loss increases with:
View Solution
Step 1: Corona loss is a phenomenon in high-voltage transmission lines where ionization of air around the conductor leads to power dissipation. Step 2: The empirical formula for corona power loss per unit length is given by: \[ P_c = f (V - V_c)^2 \frac{1}{\delta r} \] where:
- \( f \) = supply frequency,
- \( V \) = operating voltage,
- \( V_c \) = critical disruptive voltage,
- \( \delta \) = air density factor,
- \( r \) = conductor radius. Step 3: From the formula, corona loss increases when:
- Frequency (\( f \)) increases: Since \( P_c \) is directly proportional to frequency, higher supply frequency increases corona loss.
- Conductor radius (\( r \)) decreases: Smaller conductor size leads to higher electric field intensity, increasing ionization and corona effect. Step 4: Thus, the correct answer is decrease in conductor size and increase in supply frequency. Quick Tip: Corona loss can be reduced by using larger conductor diameters and operating at lower frequencies.
Which of the following matrices reveals the topology of the power system network?
View Solution
Step 1: The topology of a power system describes the interconnection of buses and branches in the network. Step 2: The bus incidence matrix (\( A \)) is a fundamental matrix in power system analysis that represents the connectivity between nodes (buses) and branches. Step 3: It is defined as: \[ A(i, j) = \begin{cases} 1, & \text{if branch } j \text{ enters bus } i
-1, & \text{if branch } j \text{ leaves bus } i
0, & \text{otherwise} \end{cases} \] Step 4: Other matrices such as:
- Primitive impedance/admittance matrices represent electrical parameters but do not describe the network topology.
- Bus admittance matrix is derived from the bus incidence matrix and does not directly reveal topology. Quick Tip: The bus incidence matrix is a key tool in power system analysis to determine network topology and formulate system equations.
Four identical alternators each rated for 20 MVA, 11 kV having a sub-transient reactance of 16% are working in parallel. The short circuit level at the bus bars is:
View Solution
Step 1: The short-circuit level (SCL) is given by: \[ SCL = \frac{\text{Total MVA Rating of Alternators}}{\text{Per Unit Sub-transient Reactance (X''})} \] Step 2: Given:
- Each alternator rating = \( 20 \) MVA
- Number of alternators = \( 4 \)
- Total MVA rating = \( 4 \times 20 = 80 \) MVA
- Sub-transient reactance (\( X'' \)) = \( 16% = 0.16 \) Step 3: Applying the formula: \[ SCL = \frac{80}{0.16} = 200 \text{ MVA} \] Step 4: Thus, the short circuit level at the bus bars is 200 MV(A) Quick Tip: The short-circuit level at bus bars is inversely proportional to the sub-transient reactance. Lower reactance results in a higher short-circuit level.
Magnetizing inrush current is rich in:
View Solution
Step 1: Magnetizing inrush current occurs in transformers when they are energized, leading to a sudden surge of current due to core saturation. Step 2: The inrush current is non-sinusoidal and contains significant harmonic components, predominantly even harmonics, particularly the 2nd harmoni(C) Step 3: The presence of 2nd harmonic is explained by:
- Flux asymmetry: When a transformer is switched on at a random point in the AC cycle, residual flux causes an asymmetrical magnetizing current.
- Nonlinear magnetization curve: The core operates in a non-linear region, generating strong even harmonics, especially the 2nd harmoni(C) Step 4: Higher-order harmonics (3rd, 5th, 7th) are present but with lower magnitude, whereas the 2nd harmonic is dominant in inrush conditions. Quick Tip: Transformer inrush current is rich in 2nd harmonic, which is used in differential protection schemes to distinguish inrush from internal faults.
Negative phase sequence current in an alternator produces:
View Solution
Step 1: Negative phase sequence (NPS) current occurs in an alternator due to unbalanced loads, leading to asymmetrical currents in the stator. Step 2: These negative sequence currents rotate in the opposite direction to the rotor’s rotation, inducing high-frequency currents in the rotor. Step 3: The impact of NPS current includes:
- Severe rotor heating due to eddy currents and I²R losses.
- Possible rotor damage if prolonged exposure occurs.
- Reduced efficiency and lifespan of the alternator. Step 4: Other options:
- Over speed (Wrong): Alternator speed is determined by the prime mover and loa(D)
- Over voltage (Wrong): Voltage variations are caused by excitation control, not NPS currents.
- Under frequency (Wrong): Frequency is controlled by grid stability, not NPS currents. Quick Tip: Negative phase sequence currents in alternators induce rotor heating and can cause serious damage, requiring protection schemes like NPS relays.
SVC is basically:
View Solution
Step 1: Static Var Compensator (SVC) is a type of Flexible AC Transmission System (FACTS) device used for voltage regulation and reactive power compensation. Step 2: SVC operates in shunt configuration, meaning it is connected in parallel to the transmission line and dynamically adjusts the reactive power. Step 3: Functionality:
- Absorbs or injects reactive power to maintain voltage stability.
- Uses thyristor-switched reactors (TSR) and thyristor-switched capacitors (TSC) for dynamic control. Step 4: Evaluation of options:
- (A) Incorrect: SVC is not a series-connected transformer-based controller.
- (B) Incorrect: SVC controls reactive power, not real power exchange.
- (C) Incorrect: SVC is a shunt-connected FACTS controller, not series. Quick Tip: Static Var Compensators (SVC) are used in power systems for voltage stabilization by dynamically adjusting reactive power in a shunt configuration.
The Impulse Response of an initially relaxed linear system is \( e^{-2t}u(t) \). To produce a response of \( te^{-2t}u(t) \), the input should be:
View Solution
Step 1: The response of a linear time-invariant (LTI) system is given by the convolution of the input \( x(t) \) with the impulse response \( h(t) \): \[ y(t) = x(t) * h(t) \] Step 2: Given:
- Impulse response: \( h(t) = e^{-2t} u(t) \)
- Desired response: \( y(t) = te^{-2t} u(t) \) Step 3: The Laplace Transform of the given impulse response is: \[ H(s) = \frac{1}{s + 2} \] And for the desired response: \[ Y(s) = \frac{1}{(s+2)^2} \] Step 4: The required input is: \[ X(s) = Y(s) / H(s) = \frac{1}{(s+2)^2} \times (s+2) = \frac{1}{s+2} \] which corresponds to: \[ x(t) = 0.5 e^{-2t} u(t) \] Quick Tip: The input required to produce a scaled and time-multiplied response in an LTI system can be determined using Laplace domain division.
The steady-state error due to unit acceleration input for a type 2 system is:
View Solution
Step 1: The steady-state error (\( e_{ss} \)) for a given system is determined using the Final Value Theorem: \[ e_{ss} = \lim_{s \to 0} s E(s) \] Step 2: For a Type 2 system, the open-loop transfer function contains two poles at the origin, meaning it has two integrators. Step 3: The steady-state error for different inputs is determined using error constants: - Position error constant \( K_p \) (for step input) - Velocity error constant \( K_v \) (for ramp input) - Acceleration error constant \( K_a \) (for parabolic input) Step 4: For a unit acceleration input (\( R(s) = \frac{1}{s^3} \)), the steady-state error is given by: \[ e_{ss} = \frac{1}{K_a} \] where \( K_a \) is the acceleration error constant. Step 5: Since a Type 2 system has two integrators, it can track acceleration inputs with finite error, given by \( \frac{1}{K_a} \). Quick Tip: For a Type 2 system, the steady-state error for an acceleration input is finite and given by \( \frac{1}{K_a} \), whereas for a Type 1 system, it would be infinite.
A system has two zeros and four poles. Then two asymptotes in the root loci plane move towards infinity along:
View Solution
Step 1: The number of root locus asymptotes is given by: \[ N_a = P - Z \] where \( P \) is the number of poles and \( Z \) is the number of zeros. Step 2: Given: - \( P = 4 \) (number of poles) - \( Z = 2 \) (number of zeros) \[ N_a = 4 - 2 = 2 \] Step 3: The angles of asymptotes are determined using the formula: \[ \theta_k = \frac{(2k + 1) 180^\circ}{N_a}, \quad k = 0, 1, \dots, (N_a - 1) \] Step 4: Substituting \( N_a = 2 \): \[ \theta_0 = \frac{(2 \times 0 + 1) 180^\circ}{2} = 90^\circ \] \[ \theta_1 = \frac{(2 \times 1 + 1) 180^\circ}{2} = -90^\circ \] Step 5: Thus, the two asymptotes move towards infinity along \( \pm 90^\circ \). Quick Tip: The angles of root locus asymptotes are calculated using the formula \( \theta_k = \frac{(2k + 1) 180^\circ}{P - Z} \). This helps in determining system stability.
A closed-loop system has the characteristic equation given by: \[ s^3 + k s^2 + (k+2) s + 3 = 0 \] For the system to be stable, the value of \( k \) is:
View Solution
Step 1: A system is stable if all the roots of the characteristic equation have negative real parts. We use Routh-Hurwitz criterion to determine stability. Step 2: Constructing the Routh array for: \[ s^3 + k s^2 + (k+2) s + 3 = 0 \] \[ \begin{array}{c|cc} s^3 & 1 & k+2
s^2 & k & 3
s^1 & \frac{k(C) - (k+2)(k)}{k} & 0
s^0 & 3 &
\end{array} \] Step 3: For stability, all first column elements must be positive. - \( k > 0 \) (ensures second row is positive) - The third-row term must be positive: \[ \frac{3k - k^2 - 2k}{k} > 0 \] \[ \frac{3k - k^2 - 2k}{k} = \frac{k - k^2}{k} > 0 \] \[ k - k^2 > 0 \] \[ k (1 - k) > 0 \] Step 4: The inequality holds for: \[ 0 < k < 1 \] Step 5: Thus, the system remains stable when \( 0 < k < 1 \). Quick Tip: The Routh-Hurwitz criterion helps determine the stability of a control system by ensuring no sign changes in the first column of the Routh array.
Loop transfer function of a feedback system is: \[ G(s)H(s) = \frac{10}{s - 2} \] Assume the Nyquist contour in the clockwise direction. Then the Nyquist plot of \( G(s) \) encircles \(-1 + j0\):
View Solution
Step 1: The Nyquist criterion determines the stability of a closed-loop system based on the open-loop transfer function \( G(s)H(s) \). Step 2: The given transfer function: \[ G(s)H(s) = \frac{10}{s - 2} \] has a single pole at \( s = 2 \), which is in the right-half plane (unstable region). Step 3: The Nyquist contour encloses the entire right half-plane, and the Nyquist plot is obtained by substituting: \[ s = j\omega \] Step 4: Evaluating at \( s = j\omega \): \[ G(j\omega)H(j\omega) = \frac{10}{j\omega - 2} \] - At high frequencies (\( \omega \to \infty \)), \( G(j\omega) \to 0 \).
- At low frequencies (\( \omega \to 0 \)), \( G(j\omega) \to -5 \).
- As \( s \) moves along the Nyquist contour, the plot encircles -1 in a clockwise direction exactly once. Step 5: According to Nyquist stability criterion: \[ N = P - Z \] where:
- \( N \) = Number of encirclements of \(-1\) (to be determined)
- \( P = 1 \) (one pole in the right-half plane)
- \( Z = 0 \) (no closed-loop unstable poles for stability) Step 6: Since \( N = 1 \), the Nyquist plot encircles \(-1\) once in a clockwise direction. Quick Tip: The Nyquist plot helps determine system stability by analyzing encirclements of \(-1\). If the plot encircles \(-1\) in clockwise direction \( P \) times, the system is unstable.
The transfer function of a first-order controller is given as: \[ G_c(s) = K \frac{(s + a)}{(s + b)} \] where \( K, a, b \) are positive real numbers. The condition for this controller to act as a phase lag compensator is:
View Solution
Step 1: A phase lag compensator is used to improve the steady-state accuracy of a control system while reducing bandwidth and increasing stability margins. Step 2: The general form of a phase lag compensator is: \[ G_c(s) = K \frac{(s + a)}{(s + b)} \] where: - \( a \) represents the zero of the compensator. - \( b \) represents the pole of the compensator. Step 3: Phase lag compensators are characterized by:
- A pole (\( b \)) closer to the origin than the zero (\( a \)).
- This ensures that at lower frequencies, the compensator reduces the phase angle, introducing a negative phase shift. Step 4: The condition for a phase lag compensator is: \[ a < b \] which ensures that the pole is dominant and the system experiences phase lag. Quick Tip: For a compensator to introduce phase lag, the pole must be closer to the origin than the zero, ensuring \( a < b \).
The state variable description of a system is: \[ \dot{X} = AX + BU, \quad A = \begin{bmatrix} 0 & 3
3 & 0 \end{bmatrix} \] The poles of the system are located at:
View Solution
Step 1: The poles of the system are given by the eigenvalues of the matrix \( A \). Step 2: The characteristic equation is found by solving: \[ \det(A - \lambda I) = 0 \] Step 3: Substituting \( A = \begin{bmatrix} 0 & 3
3 & 0 \end{bmatrix} \): \[ \begin{vmatrix} 0 - \lambda & 3
3 & 0 - \lambda \end{vmatrix} = 0 \] Step 4: Computing the determinant: \[ (-\lambda)(-\lambda) - (3 \times 3) = \lambda^2 - 9 = 0 \] Step 5: Solving for \( \lambda \): \[ \lambda^2 = 9 \] \[ \lambda = \pm j3 \] Step 6: The poles of the system are at \( s = \pm j3 \). Quick Tip: The eigenvalues of the state matrix \( A \) determine the system poles. If eigenvalues are purely imaginary \( \pm j\omega \), the system exhibits oscillatory behavior.
In a single-phase semi-converter without a freewheeling diode, for discontinuous conduction and extinction angle \( \beta > \pi \), each SCR conducts for the period:
View Solution
Step 1: In a single-phase semi-converter, SCRs (Silicon-Controlled Rectifiers) are used for controlled rectification of AC power. Step 2: The firing angle \( \alpha \) determines when the SCR turns on, and the extinction angle \( \beta \) defines when it turns off. Step 3: Without a freewheeling diode, the SCR remains conducting until the current naturally falls to zero, which occurs at the extinction angle \( \beta \). Step 4: The total conduction period of each SCR is given by: \[ \text{Conduction time} = \beta - \alpha \] Step 5: Evaluating options:
- (A) \( \pi - \alpha \) (Incorrect): This is valid only for continuous conduction mode.
- (B) \( \beta - \alpha \) (Correct): This correctly accounts for the discontinuous conduction mode where conduction stops at \( \beta \).
- (C) \( \alpha \) (Incorrect): The firing angle does not directly indicate conduction time.
- (D) \( \beta \) (Incorrect): The conduction duration is relative to \( \alpha \), not just \( \beta \). Quick Tip: In discontinuous conduction mode of a semi-converter, the conduction period of an SCR is given by \( \beta - \alpha \), where \( \alpha \) is the firing angle and \( \beta \) is the extinction angle.
For a single-phase full-wave uncontrolled rectifier with a purely \( R \) load, the form factor is:
View Solution
Step 1: The form factor (\( FF \)) of a rectifier is defined as the ratio of the RMS value of output voltage to the average output voltage: \[ FF = \frac{V_{\text{rms}}}{V_{\text{avg}}} \] Step 2: For a single-phase full-wave uncontrolled rectifier with a purely resistive (\( R \)) load: - The RMS value of the output voltage is: \[ V_{\text{rms}} = \frac{V_m}{\sqrt{2}} \] - The average output voltage is: \[ V_{\text{avg}} = \frac{2V_m}{\pi} \] Step 3: Calculating the form factor: \[ FF = \frac{V_{\text{rms}}}{V_{\text{avg}}} = \frac{\frac{V_m}{\sqrt{2}}}{\frac{2V_m}{\pi}} \] \[ FF = \frac{V_m}{\sqrt{2}} \times \frac{\pi}{2V_m} = \frac{\pi}{2\sqrt{2}} \] \[ FF = \frac{2\sqrt{2}}{\pi} \] Step 4: Thus, the correct form factor is \( \frac{2\sqrt{2}}{\pi} \). Quick Tip: The form factor helps analyze rectifier performance by comparing RMS and average voltage values. For a full-wave rectifier, \( FF = \frac{2\sqrt{2}}{\pi} \).
A single-phase inverter has a square wave output voltage. The percentage of the fifth harmonic component in relation to the fundamental component is:
View Solution
Step 1: The Fourier series representation of a square wave consists of only odd harmonics given by: \[ V_n = \frac{4V_m}{n\pi}, \quad n = 1, 3, 5, 7, \dots \] where:
- \( V_n \) is the \( n \)th harmonic component.
- \( V_m \) is the peak value of the square wave.
- \( n \) is the harmonic order. Step 2: The fundamental component (\( n = 1 \)) is: \[ V_1 = \frac{4V_m}{\pi} \] Step 3: The fifth harmonic component (\( n = 5 \)) is: \[ V_5 = \frac{4V_m}{5\pi} \] Step 4: The percentage of the fifth harmonic relative to the fundamental is: \[ \frac{V_5}{V_1} \times 100 = \frac{\frac{4V_m}{5\pi}}{\frac{4V_m}{\pi}} \times 100 \] \[ = \frac{1}{5} \times 100 = 20% \] Step 5: Thus, the percentage of the fifth harmonic component relative to the fundamental is 20%. Quick Tip: For a square wave output, harmonics follow an inverse relation with their order: \( V_n = \frac{4V_m}{n\pi} \). The fifth harmonic is always 20% of the fundamental.
The RMS output voltage at fundamental frequency of a single-phase, full-bridge inverter with input voltage of 48V DC, feeding a load of 2.4 \( \Omega \) is:
View Solution
Step 1: The output voltage waveform of a single-phase full-bridge inverter is a square wave. Step 2: The fundamental RMS output voltage (\( V_{1,\text{rms}} \)) is given by: \[ V_{1,\text{rms}} = \frac{4V_{\text{dc}}}{\sqrt{2} \pi} \] where:
- \( V_{\text{dc}} = 48V \) (DC input voltage)
- \( \pi \) is due to the Fourier series expansion of a square wave. Step 3: Substituting the given values: \[ V_{1,\text{rms}} = \frac{4 \times 48}{\sqrt{2} \pi} \text{ V} \] Step 4: Thus, the correct answer is \( \frac{4 \times 48}{\sqrt{2} \pi} \) V. Quick Tip: The fundamental RMS output voltage of a single-phase full-bridge inverter is given by \( V_{1,\text{rms}} = \frac{4V_{\text{dc}}}{\sqrt{2} \pi} \), derived from the Fourier series of a square wave.
When the MOSFET is in the ON state, the channel of the device behaves like:
View Solution
Step 1: A MOSFET (Metal-Oxide-Semiconductor Field-Effect Transistor) operates in three regions:
- Cutoff Region: Acts as an open switch.
- Linear (Ohmic) Region: Acts as a variable resistor.
- Saturation Region: Acts as a current source. Step 2: When the MOSFET is fully ON, it is in the linear region, where the channel behaves like a low-resistance path. Step 3: The ON-state resistance is called Drain-to-Source ON Resistance (\( R_{DS(on)} \)), which is relatively constant and depends on:
- Channel doping concentration.
- Gate drive voltage. Step 4: Evaluating options:
- (A) Constant resistance (Correct): The MOSFET channel exhibits ohmic behavior.
- (B) Inductance (Incorrect): The channel does not exhibit significant inductive effects.
- (C) Capacitance (Incorrect): Capacitances exist across terminals but do not dominate ON-state behavior.
- (D) Resistance and Inductance (Incorrect): The MOSFET channel primarily behaves as a resistor in the ON-state. Quick Tip: When fully ON, a MOSFET behaves like a constant low-resistance path (\( R_{DS(on)} \)), making it ideal for switching applications.
The duty cycle value of a buck converter when the switching frequency is 250 kHz and the ON time is 2\( \mu s \) is:
View Solution
Step 1: The duty cycle (\( D \)) of a buck converter is given by: \[ D = \frac{T_{\text{ON}}}{T_{\text{total}}} \] where:
- \( T_{\text{ON}} = 2 \mu s \) (ON time of the switch),
- \( T_{\text{total}} = \frac{1}{f_s} \) (Total switching period),
- \( f_s = 250 \) kHz (Switching frequency). Step 2: Calculating the total switching period: \[ T_{\text{total}} = \frac{1}{250 \times 10^3} = 4 \mu s \] Step 3: Computing the duty cycle: \[ D = \frac{2 \mu s}{4 \mu s} = 0.2 \] Step 4: Thus, the correct answer is 0.2. Quick Tip: The duty cycle of a buck converter is given by \( D = \frac{T_{\text{ON}}}{T_{\text{total}}} \). Increasing the duty cycle increases the output voltage.
Which load torque will be used in regenerative braking?
View Solution
Step 1: Regenerative braking occurs when the motor acts as a generator, converting kinetic energy into electrical energy, which is fed back into the power supply. Step 2: This braking method is effective when the load torque opposes the direction of rotation, allowing the generated energy to be recovere(D) Step 3: Passive load torque is the correct type of torque for regenerative braking because:
- It continuously applies torque in the opposite direction to rotation.
- Common in elevators, hoists, and downhill conveyors, where stored energy can be converted into electricity. Step 4: Evaluating other options:
- Fan hype load torque (Incorrect): Likely a typographical error; fans generally do not contribute to regenerative braking.
- Frictional load torque (Incorrect): Frictional torque dissipates energy as heat and does not contribute to regenerative braking.
- Archive load torque (Incorrect): This term does not apply to braking systems. Quick Tip: Regenerative braking is most effective with passive load torque, which allows energy recovery by converting mechanical motion into electrical power.
Assuming 3 MHz clock frequency, the execution time taken by the delay subroutine is:
Delay : MVI C, 9Ah
Loop : DCR C
\hspace{1.3cmJNZ Loop
\hspace{01.3cmtext{RET
View Solution
Step 1: The given program executes the following instructions: - MVI C, 9Ah (Immediate Load) → 1 Machine Cycle - DCR C (Decrement C Register) → 1 Machine Cycle per iteration - JNZ Loop (Jump if Not Zero) → 2 Machine Cycles per iteration (except last loop) Step 2: The loop executes until C = 0. Given that: \[ C = 9Ah = 154_{10} \] Step 3: The total cycles for the loop: \[ (154 - 1) \times (1 + 2) + 1 = 153 \times 3 + 1 = 460 \text{ cycles} \] Step 4: The system clock is 3 MHz, so: \[ \text{Time per cycle} = \frac{1}{3 \times 10^6} = 0.333 \mu s \] Step 5: The total execution time: \[ T = 460 \times 0.333 \mu s = 0.723 \text{ msec} \] Step 6: Thus, the correct answer is 0.723 mse(C) Quick Tip: The total delay of a loop in microprocessor programming is calculated using Clock Cycles × Time per Cycle. Always consider the number of machine cycles per instruction.
The output of the following program is: \begin{verbatim LXI H, 1234h MVI C, 05h MVI B, 67h DCR C DAD B SHLD Result HLT \end{verbatim
View Solution
Step 1: Understanding the instructions:
- LXI H, 1234h → Load register pair HL with 1234h.
- MVI C, 05h → Load C with 05h.
- MVI B, 67h → Load B with 67h.
- DCR C → Decrement C to 04h.
- DAD B → Add HL = HL + B(C) Step 2: Calculating the final value of HL:
- Initially: \( HL = 1234h \), \( BC = 6704h \).
- DAD B: \[ HL = 1234h + 6704h = 7938h \] - SHLD Result → Store HL at memory location Result.
- HLT → Halt execution. Step 3: The correct output stored in memory is 7938h. Quick Tip: The DAD B instruction in 8085 adds the contents of register pair BC to HL without affecting flags except the carry flag.
On execution of the program segment: \begin{verbatim MVI A, 0Ah SIM \end{verbatim
View Solution
Step 1: The SIM (Set Interrupt Mask) instruction in 8085 is used to control maskable interrupts using Accumulator (A) Step 2: The bit pattern of Accumulator A determines which interrupts are enabled or disabled: \[ \text{D7} \quad \text{D6} \quad \text{D5} \quad \text{D4} \quad \text{D3} \quad \text{D2} \quad \text{D1} \quad \text{D0} \] \[ \quad SOD \quad \quad \text{X} \quad \quad \text{R7.5} \quad \quad MSE \quad \quad M75 \quad \quad M65 \quad \quad M55 \] where:
- D3 (MSE - Mask Set Enable) enables the mask settings.
- D2 (M75) masks RST 7.5.
- D1 (M65) masks RST 6.5.
- D0 (M55) masks RST 5.5. Step 3: The given instruction: \[ MVI A, 0Ah = 0000\ 1010_2 \] - \( D3 = 1 \) → Masking enable(D)
- \( D1 = 0 \) → RST 6.5 not masked (enabled).
- \( D0 = 1 \) → RST 5.5 masked (disabled). Step 4: Evaluating options:
- (A) Incorrect: RST 6.5 is enable(D)
- (B) Incorrect: RST 7.5 is not affecte(D)
- (C) Correct: RST 5.5 is disabled, but other interrupts are enable(D)
- (D) Incorrect: RST 6.5 is not disable(D) Quick Tip: The SIM instruction in 8085 controls interrupt masking. When MVI A, 0Ah is executed, only RST 5.5 is disabled, while other interrupts remain enable(D)
The 8051 program segment, which performs ‘software polling’ to check if the Timer-0 counting has completed, is:
View Solution
Step 1: In the 8051 microcontroller, Timer Flags (TF0, TF1) are set when the respective timers (Timer-0, Timer-1) complete their counting. Step 2: Software polling involves continuously checking a flag until it becomes set, indicating completion. Step 3: The correct instruction format for polling is: \[ \text{JB TFx, address} \] where:
- JB (Jump if Bit is Set) checks if TF0 = 1.
- If TF0 is set, it jumps to 0FEh (execution continues).
- If TF0 is not set, it keeps polling. Step 4: Evaluating options:
- (A) JNB TF0, 0FEh (Incorrect): Jumps if TF0 is NOT set, which is opposite to polling behavior.
- (B) JB TF0, 0FEh (Correct): Checks if TF0 = 1, making it the correct software polling approach.
- (C) JB TF1, 0FEh (Incorrect): Monitors Timer-1 flag instead of Timer-0.
- (D) JNB TF1, 0FEh (Incorrect): Jumps if TF1 is NOT set, which is not polling. Quick Tip: In 8051 software polling, the correct instruction is JB TFx, address, where the program loops until the Timer Flag (TF0 or TF1) is set.
The output of the following 8051 Assembly code is: \begin{verbatim MOV A, #10 MOV 01H, A MOV A, #20 MOV @R1, A \end{verbatim
View Solution
Step 1: Analyzing each instruction: 1. MOV A, \#10 → Load immediate value 10H into Accumulator (A). \[ A = 10H \] 2. MOV 01H, A → Store the value of A (10H) into memory address 01H. \[ \text{Memory[01H]} = 10H \] 3. MOV A, \#20 → Load immediate value 20H into Accumulator (A). \[ A = 20H \] 4. MOV @R1, A → Store the value of A (20H) at the memory location pointed by Register R1. Step 2: The value of R1 is not explicitly initialized in the code, but if we assume R1 = 01H (default assumption in 8051), then: \[ \text{Memory[01H]} = 20H \] Step 3: Evaluating options:
- (A) A = 10 (Incorrect): A is 20H at the en(D)
- (B) [01] = 20 (Correct): Memory location 01H holds 20H.
- (C) [10] = 20 (Incorrect): Memory 10H was never modifie(D)
- (D) [20] = 10 (Incorrect): Memory 20H is not involve(D) Quick Tip: In 8051 Assembly, indirect addressing with @R1 stores data in the memory location pointed by R1. If R1 = 01H, then MOV @R1, A stores A's value at memory[01H].
What is the operation carried out by the 8051 instruction: ‘SETB 0D3’?
View Solution
Step 1: The instruction SETB 0D3 means "Set Bit D3 in the Special Function Register (SFR)". Step 2: In 8051 microcontrollers, bit D3 (PCON.7 - SMOD bit) in the Power Control Register (PCON) is responsible for doubling the baud rate of the serial communication. Step 3: The PCON Register (Power Control Register) structure is: \[ \begin{array}{|c|c|c|c|c|c|c|c|c|} \hline SMOD & - & - & - & GF1 & GF0 & PD & IDL
\hline D7 & D6 & D5 & D4 & D3 & D2 & D1 & D0
\hline \end{array} \] - D7 (SMOD bit) = 1 → Baud rate doubles. Step 4: Evaluating options:
- (A) Incorrect: Interrupt disable is controlled by the IE (Interrupt Enable) register.
- (B) Correct: Setting SMOD (D3 in PCON) doubles the baud rate.
- (C) Incorrect: Bank switching is done via the PSW (Program Status Word) register.
- (D) Incorrect: Timer-0 mode selection is done via the TMOD register. Quick Tip: The SMOD bit in the PCON register controls the baud rate. Using SETB PCON.7 doubles the baud rate for serial communication in 8051 microcontrollers.
If all the poles of \( H(z) \) are outside the unit circle, then the system is said to be:
View Solution
Step 1: In Discrete-Time Systems, stability and causality are analyzed using the Z-transform. Step 2: The system is BIBO (Bounded-Input Bounded-Output) stable if: \[ \sum_{n=-\infty}^{\infty} |h(n)| < \infty \] Step 3: Stability Condition in the Z-Domain:
- The system is BIBO stable if all poles lie inside the unit circle (\( |z| < 1 \)).
- If all poles are outside the unit circle, the system is not BIBO stable. Step 4: Causality Condition:
- A system is causal if its Region of Convergence (ROC) is outside the outermost pole.
- However, if all poles are outside the unit circle, the ROC is not valid for causality in practical systems. Step 5: Evaluating options:
- (A) Incorrect: The system is not necessarily causal.
- (B) Incorrect: The system is not BIBO stable.
- (C) Incorrect: The system is neither BIBO stable nor causal.
- (D) Correct: Since the system is neither BIBO stable nor causal, the correct choice is None of the above. Quick Tip: For a discrete-time system to be BIBO stable, all poles must be inside the unit circle. If all poles are outside, the system is neither BIBO stable nor causal.
Which of the following is true regarding the number of computations required to compute an \( N \)-point DFT?
View Solution
Step 1: The Discrete Fourier Transform (DFT) is computed using the formula: \[ X(k) = \sum_{n=0}^{N-1} x(n) W_N^{kn}, \quad k = 0,1,\dots,N-1 \] where \( W_N = e^{-j2\pi/N} \) is the twiddle factor. Step 2: Computational Complexity:
- Each output \( X(k) \) requires summing over \( N \) terms.
- Each term involves one complex multiplication.
- Since there are \( N \) outputs, the total number of complex multiplications is: \[ N^2 \] - Each summation involves \( (N-1) \) complex additions.
- Since there are \( N \) outputs, the total number of complex additions is: \[ N(N-1) \] Step 3: Evaluating options:
- (A) Correct: \( N^2 \) complex multiplications and \( N(N-1) \) complex additions.
- (B) Incorrect: Complex additions and multiplications are swappe(D)
- (C) Incorrect: Addition count is incorrectly given as \( N(N+1) \).
- (D) Incorrect: Both operations are incorrectly counte(D) Quick Tip: For an \( N \)-point DFT, the computational complexity is: - \( N^2 \) complex multiplications. - \( N(N-1) \) complex additions. For FFT, this is reduced to \( O(N \log N) \).
Which of the following justifies the linearity property of the \( z \)-transform? \( x(n) \leftrightarrow X(z) \).
View Solution
Step 1: The linearity property of the \( z \)-transform states: \[ \mathcal{Z} \{ a x(n) + b y(n) \} = a X(z) + b Y(z) \] where: - \( x(n) \leftrightarrow X(z) \), - \( y(n) \leftrightarrow Y(z) \), - \( a \) and \( b \) are constants. Step 2: Evaluating the given options:
- (A) Incorrect: The \( z \)-transform of a sum is the sum of individual transforms, not their product.
- (B) Correct: \( x(n) + y(n) \) transforms to \( X(z) + Y(z) \), satisfying linearity.
- (C) Incorrect: The product of signals in time does not result in an addition of their transforms.
- (D) Incorrect: The product of signals in time corresponds to convolution in the \( z \)-domain, not multiplication. Quick Tip: The linearity property of the \( z \)-transform states that \( \mathcal{Z} \{ x(n) + y(n) \} = X(z) + Y(z) \), meaning the transformation preserves addition.
What is the width of the main lobe of the frequency response of a rectangular window of length \( M-1 \)?
View Solution
Step 1: The frequency response of a rectangular window is given by the sinc function: \[ W(f) = \frac{\sin(\pi M f)}{\sin(\pi f)} \] where \( M \) is the window length. Step 2: The main lobe width of the sinc function is determined by the first zero crossings, which occur at: \[ f = \pm \frac{1}{M} \] Step 3: The total width of the main lobe in the frequency domain is: \[ \Delta \omega = \frac{4\pi}{M} \] Step 4: Evaluating options:
- (A) Incorrect: \( \frac{\pi}{M} \) is too narrow.
- (B) Incorrect: \( \frac{2\pi}{M} \) does not match the main lobe width.
- (C) Correct: \( \frac{4\pi}{M} \) matches the correct main lobe width.
- (D) Incorrect: \( \frac{8\pi}{M} \) is too wide. Quick Tip: For a rectangular window, the main lobe width of the frequency response is \( \frac{4\pi}{M} \). The wider the window, the narrower the main lobe.
With reference to the Fast Fourier Transform, if \( W_4^1 = W_x^2 \), then what is the value of \( x \)?
View Solution
Step 1: The twiddle factor in the FFT is defined as: \[ W_N = e^{-j\frac{2\pi}{N}} \] where \( W_N^k \) represents the \( k \)th power of the twiddle factor for an \( N \)-point FFT. Step 2: Given: \[ W_4^1 = W_x^2 \] Substituting the definition of the twiddle factor: \[ e^{-j\frac{2\pi}{4} \times 1} = e^{-j\frac{2\pi}{x} \times 2} \] Step 3: Simplifying: \[ e^{-j\frac{2\pi}{4}} = e^{-j\frac{4\pi}{x}} \] Equating exponents: \[ \frac{2\pi}{4} = \frac{4\pi}{x} \] \[ \frac{\pi}{2} = \frac{4\pi}{x} \] Solving for \( x \): \[ x = 4 \] Step 4: Evaluating options:
- (A) Incorrect: \( x = 2 \) does not satisfy the equation.
- (B) Correct: \( x = 4 \) is the correct answer.
- (C) Incorrect: \( x = 8 \) does not match.
- (D) Incorrect: \( x = 16 \) is incorrect. Quick Tip: The twiddle factor in FFT is given by \( W_N = e^{-j\frac{2\pi}{N}} \). When equating twiddle factors, always compare exponents carefully.
Which of the following defines the FIR filter for length \( M \), input \( x(n) \), and output \( y(n) \)?
View Solution
Step 1: A Finite Impulse Response (FIR) filter is defined by the convolution sum: \[ y(n) = \sum_{k=0}^{M-1} b_k x(n-k) \] where:
- \( x(n) \) is the input signal,
- \( y(n) \) is the output signal,
- \( b_k \) are the filter coefficients,
- \( M \) is the filter length. Step 2: FIR filters are non-recursive, meaning they only depend on current and past input values. Step 3: Evaluating the given options:
- (A) Incorrect: The upper limit should be \( M-1 \), not \( M \).
- (B) Incorrect: FIR filters use past values, not future values (\( x(n+k) \)).
- (C) Correct: Matches the standard FIR filter definition.
- (D) Incorrect: Uses \( x(n+k) \), which does not define a standard FIR filter.
Quick Tip: A FIR filter is implemented using convolution: \( y(n) = \sum_{k=0}^{M-1} b_k x(n-k) \). It depends only on past and current inputs.
Surge impedance of a lossless transmission line is (if \( L \) = inductance/m and \( C \) = capacitance/m):
View Solution
Step 1: The characteristic impedance (or surge impedance) of a lossless transmission line is given by: \[ Z_0 = \sqrt{\frac{L}{C}} \] where:
- \( L \) is the inductance per unit length (H/m),
- \( C \) is the capacitance per unit length (F/m). Step 2: This formula is derived from the transmission line equation, considering a lossless line where resistance (\( R \)) and conductance (\( G \)) are negligible. Step 3: Evaluating options:
- (A) Incorrect: The correct formula has \( L \) in the numerator, not \( C \).
- (B) Correct: \( \sqrt{\frac{L}{C}} \) is the correct surge impedance expression.
- (C) Incorrect: \( \frac{1}{\sqrt{LC}} \) is incorrect.
- (D) Incorrect: \( \sqrt{LC} \) does not represent surge impedance. Quick Tip: The surge impedance of a lossless transmission line is \( Z_0 = \sqrt{\frac{L}{C}} \). It plays a crucial role in impedance matching and wave propagation.
Time lag for breakdown is:
View Solution
Step 1: Definition of Time Lag for Breakdown The time lag for breakdown in gaseous insulation refers to the delay between the application of a voltage higher than the breakdown voltage and the actual occurrence of breakdown. Step 2: Breakdown Time Lag Components
The time lag consists of:
- Statistical time lag: Time required for the formation of a free electron to initiate breakdown.
- Formative time lag: Time taken for avalanche multiplication to result in complete breakdown.
Step 3: Evaluating options:
- (A) Incorrect: Breakdown occurs in gases under different conditions, not necessarily under pulse application.
- (B) Incorrect: Voltage rise is independent of the breakdown time lag.
- (C) Correct: Time lag is defined as the delay between applied voltage and breakdown occurrence.
- (D) Incorrect: Ionization is a part of the process but does not fully define the time lag. Quick Tip: The time lag for breakdown is the delay between voltage application and the actual breakdown, influenced by statistical and formative factors.
In impulse testing of transformers, fault location is usually done by:
View Solution
Step 1: Impulse Testing of Transformers Impulse testing is conducted to assess the insulation strength of transformers under transient voltage conditions, simulating lightning surges or switching surges. Step 2: Fault Detection Methods During impulse testing, faults are identified using oscillographic analysis of current and voltage waveforms. - Neutral Current Oscillogram:
- Used to detect internal insulation faults in transformers.
- A sudden change in the waveform indicates insulation failure. - Chopped Wave Oscillogram:
- Helps in overvoltage withstand testing but is not primarily used for fault location. Step 3: Evaluating options:
- (A) Correct: Neutral current oscillogram effectively detects faults in impulse testing.
- (B) Incorrect: Chopped wave tests insulation performance but not specific fault locations.
- (C) Incorrect: Noise or smoke observation is not a precise fault detection metho(D)
- (D) Incorrect: Scanning is not a standard practice in transformer impulse testing. Quick Tip: In transformer impulse testing, faults are located using the neutral current oscillogram, which shows waveform distortions indicating insulation failure.
Optimum number of stages for Cockcroft-Walton voltage multiplier circuit are:
View Solution
Step 1: Cockcroft-Walton Voltage Multiplier The Cockcroft-Walton (CW) multiplier is a circuit used to generate high DC voltages from an AC input using stages of capacitors and diodes. Step 2: Optimum Number of Stages The voltage drop in a CW multiplier increases with the number of stages, which affects efficiency. The optimum number of stages (\( n \)) is given by: \[ n_{\text{opt}} = \sqrt{\frac{V_{\max}}{I f C}} \] where:
- \( V_{\max} \) = Maximum output voltage,
- \( I \) = Load current,
- \( f \) = Frequency of AC supply,
- \( C \) = Capacitance of the multiplier capacitors. Step 3: Evaluating options:
- (A) Correct: Matches the standard equation for optimal stage number.
- (B) Incorrect: Inverted ratio, incorrect formul(A)
- (C) Incorrect: Incorrect placement of terms.
- (D) Incorrect: Incorrect expression for capacitance impact. Quick Tip: The optimum number of stages in a Cockcroft-Walton voltage multiplier is given by: \[ n_{\text{opt}} = \sqrt{\frac{V_{\max}}{I f C}} \] Balancing capacitance, frequency, and load current helps optimize voltage gain efficiency.
The most important test to assert the proper functions of a surge diverter is:
View Solution
Step 1: Surge Diverter (Lightning Arrester) Testing A surge diverter (lightning arrester) is used to protect electrical equipment from high transient voltages caused by lightning or switching surges. Step 2: Impulse Current Test The most critical test for verifying a surge diverter's performance is the impulse current test, which ensures:
- The diverter can withstand high impulse currents from lightning strikes.
- It operates correctly to divert excess voltage safely. Step 3: Evaluating options:
- (A) Incorrect: Impulse withstand tests evaluate insulation strength but are not the primary test for surge diverters.
- (B) Incorrect: Spark over and residual voltage tests measure breakdown voltage but do not fully validate operational performance.
- (C) Correct: The impulse current test is the primary test ensuring proper function.
- (D) Incorrect: Pollution tests check environmental resistance but are not the most crucial test. Quick Tip: The impulse current test is the most important test for a surge diverter, ensuring it can handle lightning and switching surges effectively.
An R-C voltage divider has an HV arm capacitance, \( C_1 = 600 \) pF, resistance \( R = 400 \) \( \Omega \), and equivalent ground capacitance \( C_g = 240 \) pF. The effective time constant of the divider in nanoseconds is:
View Solution
Step 1: Effective Capacitance Calculation The total effective capacitance for an R-C voltage divider is given by: \[ C_{\text{eff}} = \frac{C_1 C_g}{C_1 + C_g} \] Substituting given values: \[ C_{\text{eff}} = \frac{(600 \times 240)}{(600 + 240)} \text{ pF} \] \[ C_{\text{eff}} = \frac{144000}{840} = 171.43 \text{ pF} \] Step 2: Time Constant Calculation The time constant \( \tau \) is given by: \[ \tau = R \cdot C_{\text{eff}} \] \[ \tau = 400 \times 171.43 \text{ ps} \] \[ \tau = 68571.43 \text{ ps} = 68.57 \text{ ns} \approx 67 \text{ ns} \] Step 3: Evaluating options:
- (A) Incorrect: 32 ns is too low.
- (B) Incorrect: 100 ns is an overestimation.
- (C) Correct: \( 67 \) ns matches the computed result.
- (D) Incorrect: 25 ns is too low. Quick Tip: The effective time constant of an R-C voltage divider is computed using: \[ \tau = R \cdot \frac{C_1 C_g}{C_1 + C_g} \] For accurate impulse response analysis in high-voltage systems.
Electric traction uses a power supply of:
View Solution
Step 1: Electric Traction Power System Electric railway traction systems predominantly use single-phase 25 kV AC at 50 Hz, which is the global standard for high-speed and mainline railways. Step 2: Why 25 kV AC?
- High Voltage (25 kV) Benefits: Reduces current and minimizes power losses.
- AC vs. DC: AC systems allow for efficient power transmission over long distances.
- 50Hz Frequency: Compatible with most power grids worldwide. Step 3: Evaluating options:
- (A) Correct: 25 kV AC at 50 Hz is the standard for electric traction.
- (B) Incorrect: DC traction is used in older urban metro systems, but high-speed rail uses A(C)
- (C) Incorrect: 50 kV is not a standard traction voltage.
- (D) Incorrect: 50 kV DC is not used in railway traction. Quick Tip: Most modern railway traction systems use single-phase 25 kV AC at 50 Hz, ensuring efficient transmission and reduced power losses.
Filament lamps operate normally at a power factor of:
View Solution
Step 1: Understanding Power Factor The power factor of an electrical load is defined as: \[ \text{Power Factor} = \cos \theta = \frac{\text{Real Power}}{\text{Apparent Power}} \] Step 2: Power Factor of Filament Lamps
- Filament lamps (incandescent lamps) consist of a resistive heating element (usually tungsten).
- In a purely resistive circuit, voltage and current are in phase, leading to a unity power factor (\( \cos \theta = 1 \)). Step 3: Evaluating options:
- (A) Incorrect: A lagging power factor is associated with inductive loads (motors, transformers).
- (B) Incorrect: A leading power factor is associated with capacitive loads.
- (C) Incorrect: A zero power factor is observed in circuits dominated by pure reactance (inductors or capacitors).
- (D) Correct: Filament lamps have a unity power factor as they behave like a pure resistive loa(D) Quick Tip: Filament lamps operate at a unity power factor (\( PF = 1 \)), since they behave like a pure resistive load with no inductive or capacitive effects.
Candela is the unit of:
View Solution
Step 1: Understanding Candela (cd) The candela (cd) is the SI unit of luminous intensity, which represents the amount of light emitted by a source in a particular direction. Step 2: Differentiating Related Terms
- Luminous flux (lumen, lm): Measures the total visible light output from a source.
- Luminous intensity (candela, cd): Measures the light emitted per unit solid angle in a specific direction.
- Light: A general term, not a specific unit.
- Brightness: A subjective measure of perceived light, not an SI unit. Step 3: Evaluating options:
- (A) Incorrect: Luminous flux is measured in lumens (lm), not candel(A)
- (B) Correct: Luminous intensity is measured in candela (cd).
- (C) Incorrect: Light is a general term, not a physical unit.
- (D) Incorrect: Brightness is a subjective perception, not an SI unit. Quick Tip: The candela (cd) is the SI unit of luminous intensity, which measures the visible light emitted in a particular direction.
A slab of insulating material 130 cm\(^2\) in area and 1 cm thick is to be heated by dielectric heating. The power required is 380 W at 30 MHz. The material has a relative permittivity of 5 and a power factor of 0.05. Determine the necessary voltage.
View Solution
Step 1: Dielectric Heating Power Formula The power dissipated in dielectric heating is given by: \[ P = \pi f V^2 C \tan \delta \] where:
- \( P \) = power (380 W),
- \( f \) = frequency (30 MHz = \( 30 \times 10^6 \) Hz),
- \( V \) = applied voltage (to be determined),
- \( C \) = capacitance of the dielectric material,
- \( \tan \delta \) = power factor (0.05). Step 2: Capacitance Calculation Capacitance of a parallel plate capacitor is given by: \[ C = \frac{\varepsilon_0 \varepsilon_r A}{d} \] where:
- \( \varepsilon_0 \) = permittivity of free space (\( 8.854 \times 10^{-12} \) F/m),
- \( \varepsilon_r \) = relative permittivity (5),
- \( A \) = area (\( 130 \times 10^{-4} \) m\(^2\)),
- \( d \) = thickness (0.01 m). \[ C = \frac{(8.854 \times 10^{-12} \times 5) \times (130 \times 10^{-4})}{0.01} \] \[ C = 5.75 \times 10^{-12} \text{ F} \] Step 3: Voltage Calculation Rearranging the power equation: \[ V = \sqrt{\frac{P}{\pi f C \tan \delta}} \] Substituting the values: \[ V = \sqrt{\frac{380}{\pi \times (30 \times 10^6) \times (5.75 \times 10^{-12}) \times 0.05}} \] \[ V = \sqrt{\frac{380}{\pi \times 8.625 \times 10^{-5}}} \] \[ V = \sqrt{1.4 \times 10^7} \] \[ V \approx 837 V \] Step 4: Evaluating options: - (A) Incorrect: 837 kV is an overestimation. - (B) Correct: 837 V matches the computed result. - (C) Incorrect: 652 V is too low. - (D) Incorrect: 552 V is incorrect. Quick Tip: For dielectric heating, the required voltage is calculated using: \[ V = \sqrt{\frac{P}{\pi f C \tan \delta}} \] where capacitance is derived from material properties.
Spot welding is used for:
View Solution
Step 1: What is Spot Welding?
- Spot welding is a type of resistance welding used to join two or more metal sheets together by applying pressure and electrical current.
- It is commonly used in automobile, aerospace, and sheet metal industries. Step 2: Why Thin Metal Sheets?
- Spot welding works best for thin metal sheets (typically 0.5 mm to 3 mm thick).
- It creates localized heat at the weld points using high current, causing the sheets to fuse. Step 3: Evaluating options:
- (A) Correct: Spot welding is primarily used for thin metal sheets.
- (B) Incorrect: Thick metal rods require arc welding or friction welding.
- (C) Incorrect: Thick square sections need MIG/TIG welding for strong joints.
- (D) Incorrect: Rough or irregular surfaces are difficult to weld using spot welding. Quick Tip: Spot welding is ideal for thin metal sheets in industries like automobile and electronics due to its high-speed operation and low material distortion.
Material used for solar cell is:
View Solution
Step 1: Understanding Solar Cell Materials A solar cell (photovoltaic cell) converts sunlight into electricity through the photovoltaic effect. The efficiency and performance of solar cells depend on the material use(D) Step 2: Why Silicon?
- Silicon (Si) is the most commonly used material in solar cells due to:
- Abundance: Readily available and cost-effective.
- Semiconductor Properties: Ideal for electron flow under sunlight.
- High Efficiency: Monocrystalline and polycrystalline silicon provide high energy conversion rates. Step 3: Evaluating options:
- (A) Incorrect: Germanium is a semiconductor but less efficient than silicon in solar applications.
- (B) Correct: Silicon is the standard material for solar cells.
- (C) Incorrect: Silica gel is used for moisture absorption, not in photovoltaic cells.
- (D) Incorrect: Mercury is not used in solar cell fabrication. Quick Tip: Silicon (Si) is the primary material used in solar cells due to its high efficiency, abundance, and excellent semiconductor properties.



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